Animagine XL (legacy): v2.0: animagineXLLegacy_v20.safetensors

SDXL 1.0 Standard Model
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ANIMAGINE XL 2.0

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Huggingface link: https://huggingface.co/Linaqruf/animagine-xl-2.0

Gradio Demo : https://huggingface.co/spaces/Linaqruf/Animagine-XL

Training model is very expensive, please consider to support us at : https://ko-fi.com/linaqruf

Overview

Animagine XL 2.0 is an advanced latent text-to-image diffusion model designed to create high-resolution, detailed anime images. It's fine-tuned from Stable Diffusion XL 1.0 using a high-quality anime-style image dataset. This model, an upgrade from Animagine XL 1.0, excels in capturing the diverse and distinct styles of anime art, offering improved image quality and aesthetics.

Model Details

  • Developed by: Linaqruf

  • Model type: Diffusion-based text-to-image generative model

  • Model Description: This is a model that excels in creating detailed and high-quality anime images from text descriptions. It's fine-tuned to understand and interpret a wide range of descriptive prompts, turning them into stunning visual art.

  • License: CreativeML Open RAIL++-M License

  • Finetuned from model: Stable Diffusion XL 1.0

Gradio & Colab Integration

Animagine XL is accessible via Gradio Web UI and Google Colab, offering user-friendly interfaces for image generation:

  • Gradio Web UI: Open In Spaces

  • Google Colab: Open In Colab

Further information please visit Official Repository at Huggingface